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A Nonlinear Adiabatic Theorem for Coherent States

Rémi Carles, Clotilde Fermanian Kammerer

To cite this version:

Rémi Carles, Clotilde Fermanian Kammerer. A Nonlinear Adiabatic Theorem for Coherent States.

Nonlinearity, IOP Publishing, 2011, 24 (8), pp.2143-2164. �10.1088/0951-7715/24/8/002�. �hal-

00530175�

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STATES

R´EMI CARLES AND CLOTILDE FERMANIAN-KAMMERER

Abstract. We consider the propagation of wave packets for a one-dimensional nonlinear Schr¨odinger equation with a matrix-valued potential, in the semi- classical limit. For an initial coherent state polarized along some eigenvector, we prove that the nonlinear evolution preserves the separation of modes, in a scaling such that nonlinear effects are critical (the envelope equation is non- linear). The proof relies on a fine geometric analysis of the role of spectral projectors, which is compatible with the treatment of nonlinearities. We also prove a nonlinear superposition principle for these adiabatic wave packets.

1. Introduction

We consider the semi-classical limit ε → 0 for the nonlinear Schr¨odinger equation (1.1)

 

iε∂

t

ψ

ε

+ ε

2

2 ∂

x2

ψ

ε

= V (x)ψ

ε

+ Λ | ψ

ε

|

2CN

ψ

ε

, (t, x) ∈ R × R, ψ

|t=0ε

= ψ

0ε

where Λ ∈ R. The data ψ

0ε

and the solution ψ

ε

(t) are vectors of C

N

, N > 1. The quantity | ψ

ε

|

2CN

denotes the square of the Hermitian norm in C

N

of the vector ψ

ε

. Finally, the potential V is smooth and valued in the set of N by N Hermitian matrices. Such systems appear in the modelling of Bose-Einstein condensate (see [1]

and references therein).

Definition 1.1. We say that a function f is at most quadratic if f ∈ C

(R) and for all k > 2, f

(k)

∈ L

(R).

We make the following assumptions on the potential V :

Assumption 1.2. (1) We have V (x) = D(x)+W (x) with D, W ∈ C

(R, R

N×N

), D diagonal with at most quadratic coefficients, and W symmetric and bounded as well as its derivatives, W ∈ W

∞,∞

(R).

(2) The matrix V has P distinct, at most quadratic, eigenvalues λ

1

, . . . , λ

P

and (1.2) ∃ c

0

, n

0

∈ R

+

, ∀ j 6 = k, ∀ x ∈ R , | λ

j

(x) − λ

k

(x) | > c

0

h x i

−n0

.

Under these assumptions (the first point suffices), we can prove global existence of the solution ψ

ε

for fixed ε > 0:

Lemma 1.3. If V satisfies Assumption 1.2 and ψ

0ε

∈ L

2

( R ), there exists a unique, global, solution to (1.1)

ψ

ε

∈ C R; L

2

(R)

∩ L

8loc

R; L

4

(R) .

This work was supported by the French ANR project R.A.S. (ANR-08-JCJC-0124-01).

1

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The L

2

-norm of ψ

ε

does not depend on time: k ψ

ε

(t) k

L2(R)

= k ψ

ε0

k

L2(R)

, ∀ t ∈ R.

The proof of this lemma is sketched in Appendix A.

In this nonlinear setting, the size of the initial data is crucial. As in [4], we choose to consider initial data of order 1 (in L

2

), and to introduce a dependence upon ε in the coupling constant (note that the nonlinearity is homogeneous). This leads to the equation

iε∂

t

ψ

ε

+ ε

2

2 ∂

x2

ψ

ε

= V (x)ψ

ε

+ Λε

| ψ

ε

|

2CN

ψ

ε

,

and we choose the exponent β = 3/4, which is critical for the type of initial data we want to consider (coherent state) when the potential V is scalar (see [4]). We are left with the nonlinear semi-classical Schr¨odinger equation

(1.3) iε∂

t

ψ

ε

+ ε

2

2 ∂

x2

ψ

ε

= V (x)ψ

ε

+ Λε

3/2

| ψ

ε

|

2CN

ψ

ε

; ψ

|t=0ε

= ψ

ε0

. We focus on initial data which are perturbation of wave packets

(1.4) ψ

0ε

(x) = ε

−1/4

e

0(x−x0)/ε

a

x − x

0

√ ε

χ(x) + r

0ε

(x), where the initial error satisfies

(1.5) k r

ε0

k

L2(R)

+ k xr

0ε

k

L2(R)

+ k ε∂

x

r

ε0

k

L2(R)

= O (ε

κ

) for some κ > 1 4 . The profile a belongs to the Schwartz class, a ∈ S (R; C), and the initial datum is polarized along an eigenvector χ(x) ∈ C

(R; C

N

):

V (x)χ(x) = λ

1

(x)χ(x), with | χ(x) |

CN

= 1.

Note that λ

1

is simply a notation for some eigenvalue, up to a renumbering of eigenvalues. The L

2

-norm of ψ

ε0

is independent of ε, k ψ

0ε

k

L2(R)

= k a k

L2(R)

. As pointed out above, this is equivalent to considering (1.1) with initial data of the same form (1.4), but of order ε

3/4

in L

2

(R). The evolution of such data when a is a Gaussian has been extensively studied by G. Hagedorn on the one hand, and by G. Hagedorn and A. Joye on the other hand, in the linear context Λ = 0 (see [8, 9]).

These data are also particularly interesting for numerics (see [12] and the references therein).

Because of the gap condition, the matrix V has smooth eigenvalues and eigenpro- jectors (see [11]). Besides, the gap condition (1.2) also implies that we control the growth of the eigenprojectors (see Lemma C.2). Note however than in dimension 1 (x ∈ R ), one can have smooth eigenprojectors without any gap condition. We explain this fact below and give an example of projectors that we can consider; we also illustrate why things may be more complicated in higher dimensions (d > 2).

Example 1.4. For N = 2 and x ∈ R, consider (1.6) V (x) = (ax

2

+ b)Id +

u(x) v(x) v(x) − u(x)

,

for a, b ∈ R, and u and v smooth and bounded with bounded derivatives. Such a potential satisfies Assumption 1.2. Its eigenvalues are the two functions

λ

±

(x) = ax

2

+ b ± p

u(x)

2

+ v(x)

2

.

These functions are clearly smooth outside the set of points x

0

such that u(x

0

)

2

+

v(x

0

)

2

= 0. Besides, for such points, one can renumber the modes in order to

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build smooth eigenvalues. More precisely, observe first that if u(x)

2

+ v(x)

2

= O((x − x

0

)

) close to x

0

, the functions λ

±

are smooth close to x

0

. Moreover, if u(x)

2

+ v(x)

2

= (x − x

0

)

k

f (x) with f (x

0

) 6 = 0, necessarily f (x

0

) > 0 and k = 2p, so we have

λ

±

(x) = ax

2

+ b ± | x − x

0

|

p

p f (x).

For p even these functions are again smooth. However, when p is odd, they are no longer smooth and we perform a renumbering of the eigenfunctions, observing that

x 7→ ax

2

+ b + (x − x

0

)

p

p f (x) are smooth eigenvalues of V close to x

0

.

Example 1.5. Resume the above example, with now x ∈ R

d

, d > 2. The smoothness of the eigenvalues is no longer guaranteed: suppose u(x) = x

1

and v(x) = x

2

, then the functions λ

±

are not smooth and one cannot find any renumbering which makes them smooth.

Example 1.6. For an example of a potential which satisfies 1.2, we simply choose V as in (1.6) with

c

u

u(x) = c

v

v(x) = h x i

−n0

, c

2u

+ c

2v

6 = 0.

1.1. The ansatz. We consider the classical trajectories (x(t), ξ(t)) solutions to (1.7) x(t) = ˙ ξ(t), ξ(t) = ˙ −∇ λ

1

(x(t)), x(0) = x

0

, ξ(0) = ξ

0

.

Because λ

1

is at most quadratic, the classical trajectories grow at most exponen- tially in time (see e.g. [4]):

(1.8) ∃ C > 0, | ξ(t) | + | x(t) | . e

Ct

. We denote by S the action associated with (x(t), ξ(t))

(1.9) S(t) =

Z

t 0

1

2 | ξ(s) |

2

− λ

1

(x(s))

ds.

We consider the function u = u(t, y) solution to (1.10) i∂

t

u + 1

2 ∂

y2

u = 1

2 λ

′′1

(x(t)) y

2

u + Λ | u |

2

u ; u(0, y) = a(y), and we denote by ϕ

ε

the function associated with u, x, ξ, S by:

(1.11) ϕ

ε

(t, x) = ε

−1/4

u

t, x − x(t)

√ ε

e

i(S(t)+ξ(t)(x−x(t)))/ε

.

Global existence of u and control of its derivatives and momenta are proved in [3].

More precisely, we have the following result.

Theorem 1.7 (From [3]). Suppose a ∈ S (R). There exists a unique, global so- lution u ∈ C (R; L

2

(R)) ∩ L

8loc

(R; L

4

(R)) to (1.10). In addition, for all k, p ∈ N, h y i

k

yp

u ∈ C ( R ; L

2

( R )) and

(1.12) ∀ k, p ∈ N, ∃ C > 0, ∀ t ∈ R

+

, k h y i

k

yp

u(t, · ) k

L2(R)

. e

Ct

.

In particular, note that ∂

py

u(t, · ) is in L

for all p ∈ N. These results have

consequences on ϕ

ε

. As far as the L

norm is concerned, we infer, using (1.8),

(1.13) ∀ p ∈ N , k (ε∂

x

)

p

ϕ

ε

(t) k

L

. ε

−1/4

e

Cpt

.

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We use the time-dependent eigenvectors constructed in [8] (see also [9] and [14]).

To make the notations precise, we denote by d

j

the multiplicity of the eigenvalue λ

j

, 1 6 j 6 P (note that P

16j6P

d

j

= N).

Proposition 1.8. There exists a smooth orthonormal family χ

(t, x)

16ℓ6d1

such that for all t, χ

(t, x)

16ℓ6d1

spans the eigenspace associated to λ

1

, χ

1

(0, x) = χ(x) and for m ∈ { 1, · · · , d

1

} ,

(1.14) χ

m

(t, x), ∂

t

χ

(t, x) + ξ(t)∂

x

χ

(t, x)

CN

= 0.

Moreover, for ℓ ∈ { 1, . . . , d

1

} , k, p ∈ N, there exists a constant C = C(p, k) such that

tp

xk

χ

(t, x)

CN

6 C e

Ct

h x i

(k+p)(1+n0)

,

where n

0

appears in (1.2).

Note that equation (1.14) for m = ℓ is true as soon as the eigenvector χ

is normalized and real-valued.

Equation (1.14) is often referred to as parallel transport. These time-dependent eigenvectors are commonly used in adiabatic theory and are connected with the Berry phase (see [14]). Their construction is recalled in Section 2, where the control of their growth is also established.

Notation. In the case of a single coherent state, we complete the family χ

(t, x)

16ℓ6d1

as an orthonormal basis χ

j

16j6P 16ℓ6dj

of C

N

as follows:

• χ

1

= χ

,

• For j > 2 and 1 6 ℓ 6 d

j

, χ

j

= χ

j

(x) does not depend on time,

• For j > 2, (χ

j

)

16j6dj

spans the eigenspace associated to λ

j

.

1.2. The results. We prove that there is adiabatic decoupling for the solution of (1.3) with initial data which are coherent states of the form (1.4): the solution keeps the same form and remains in the same eigenspace.

Theorem 1.9. Let a ∈ S (R) and r

ε0

satisfying (1.5). Under Assumption 1.2, con- sider ψ

ε

solution to the Cauchy problem (1.3)–(1.4), and the approximate solution ϕ

ε

given by (1.11). There exists a constant C > 0 such that the function

w

ε

(t, x) = ψ

ε

(t, x) − ϕ

ε

(t, x)χ

1

(t, x), where χ

1

is given by Proposition 1.8, satisfies

sup

|t|6Cloglog1ε

( k w

ε

(t) k

L2

+ k xw

ε

(t) k

L2

+ k ε∂

x

w

ε

(t) k

L2

) −→

ε→0

0.

This adiabatic decoupling between the modes is well-known in the linear setting and is at the basis of numerous results on semi-classical Schr¨odinger operator with matrix-valued potential in the framework of Born-Oppenheimer approximation for molecular dynamics. On this subject, the reader can consult the article of H. Spohn and S. Teufel [13] or the book of S. Teufel [14] for a review on the topic (see also [2]

for an adiabatic result in a nonlinear context and [10] for application of adiabatic theory to the obtention of resolvent estimates).

Remark 1.10. Suppose that V depends on ε with V

ε

= D + εW , where D and W are as in Assumption 1.2; this is so in the model presented in [1]. Then the above result remains true for | t | 6 C log

1ε

: we gain one logarithm. See Remark 4.1 for

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the key arguments. Also, the assumption on the initial error can be relaxed: to prove the analogue of Theorem 1.9 with an approximation in L

2

up to C log 1/ε, (1.5) can be replaced with

k r

ε0

k

L2(R)

→ 0 as ε → 0.

In contrast with the general framework of this paper, no rate is needed: the rate in (1.5) is due to the fact that we cannot use Strichartz estimates here.

It is also interesting to analyze the evolution of solution associated with data which are the superposition of two data of the studied form. We suppose

ψ

0ε

(x) = ϕ

ε1

(0, x)χ

1

(x) + ϕ

ε2

(0, x)χ

2

(x),

where both functions ϕ

ε1

and ϕ

ε2

have the form (1.11), for two eigenvectors of V , χ

1

and χ

2

, and phase space points (x

1

, ξ

1

) and (x

2

, ξ

2

). We assume

1

, x

1

, ξ

1

) 6 = (χ

2

, x

2

, ξ

2

) .

We associate with the phase space points (x

j

, ξ

j

), j ∈ { 1, 2 } the classical trajectories (x

j

(t), ξ

j

(t)), and the action S

j

(t) associated with ˜ λ

j

such that

V (x)χ

j

(x) = ˜ λ

j

(x)χ

j

(x).

Note that we may have ˜ λ

1

= ˜ λ

2

. Let us denote by χ

j

(t)

16j6P 16ℓ6dj

a time-dependent orthonormal basis of eigenvectors defined according to Proposition 2.1 (see also Proposition 1.8 above) with χ

11

(0, x) = χ

1

(x), χ

2

(x) = χ

21

(0, x) if ˜ λ

1

= ˜ λ

2

, χ

2

(x) = χ

12

(0, x) otherwise, and by ϕ

εj

the ansatz defined by (1.11). To unify the presenta- tion, we write

χ

1

= χ

11

; χ

2

=

( χ

21

if ˜ λ

1

= ˜ λ

2

, χ

12

otherwise.

Theorem 1.11. Set E

j

=

ξ

2 j

2

+ ˜ λ

j

(x

j

) for j ∈ { 1, 2 } and suppose Γ = inf

x∈R

λ ˜

1

(x) − λ ˜

2

(x) − (E

1

− E

2

) > 0.

There exists C > 0 such that the function

w

ε

(t) = ψ

ε

(t) − ϕ

ε1

χ

1

(t, x) − ϕ

ε2

χ

2

(t, x).

satisfies

sup

t6Cloglog1ε

( k w

ε

(t) k

L2

+ k xw

ε

(t) k

L2

+ k ε∂

x

w

ε

(t) k

L2

) −→

ε→0

0.

Note that if ˜ λ

1

= ˜ λ

2

, one recovers the condition E

1

6 = E

2

of [4]. The proof of Theorem 1.11 follows the same lines as in [4, Section 6]. The constant Γ controls the frequencies of time interval where trajectories cross.

Remark 1.12. In finite time, the situation is different whether ˜ λ

1

= ˜ λ

2

or not. If λ ˜

1

= ˜ λ

2

, the superposition holds in finite time without any condition on Γ; this comes from the fact that the trajectories x

1

(t) and x

2

(t) only cross on isolated points (see [4]). However, if ˜ λ

1

6 = ˜ λ

2

one may have x

1

(t) = x

2

(t) on intervals of non-empty interior: the condition Γ 6 = 0 prevents this situation from happening.

For example, if

V (x) =

cosx sinx sinx − cosx

+ v(x)Id

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with v smooth and at most quadratic, we have λ

1

(x) = v(x) − 1 and λ

2

(x) = v(x)+1:

classical trajectories for both modes, issued from the same point of the phase space, are equal.

1.3. Strategy of the proof of Theorem 1.9. The proof is more complicated than in the scalar case [4], due to the fact that the spectral projectors do not commute with the Laplace operator. From this perspective, a much finer geometric understanding is needed and we revisit [8, 7, 9, 13, 14] by adapting to our nonlinear context ideas contained therein.

Observe first that the function ϕ

ε

satisfies (1.15) iε∂

t

ϕ

ε

+ ε

2

2 ∂

x2

ϕ

ε

= T

ε

(t, x)ϕ

ε

+ Λ ε

3/2

| ϕ

ε

|

2CN

ϕ

ε

, where

T

ε

(t, x) = λ

1

(x(t)) + λ

1

(x(t))(x − x(t)) + 1

2 λ

′′

(x(t))(x − x(t))

2

.

This term corresponds to the beginning of the Taylor expansion of λ

1

about x(t).

Therefore, the function w

ε

(t, x) = ψ

ε

(t, x) − ϕ

ε

(t, x)χ

1

(t, x) satisfies w

|t=0ε

= r

ε0

and iε∂

t

w

ε

(t, x) + ε

2

2 ∂

x2

w

ε

(t, x) − V (x)w

ε

(t, x) = ε N L g

ε

(t, x) + ε L e

ε

(t, x) where

N L g

ε

= Λ ε

1/2

ϕ

ε

χ

1

+ w

ε

2CN

ε

χ

1

+ w

ε

) − | ϕ

ε

|

2

ϕ

ε

χ

1

, L e

ε

= i∂

t

χ

1

ϕ

ε

+ ε∂

x

χ

1

x

ϕ

ε

+ ε

2 ϕ

ε

x2

χ

1

+ ε

−1

1

(x) − T

ε

) ϕ

ε

χ

1

. Since ϕ

ε

is concentrated near x = x(t) at scale √

ε, we have (λ

1

(x) − T

ε

ε

= O

ε

3/2

e

Ct

in L

2

(R),

where we have used Theorem 1.7. The term L e

ε

a priori presents an O (1) con- tribution, which is an obstruction to infer that w

ε

is small by applying Gronwall Lemma. Observing that in view of the estimates on the classical flow (see (1.8))

ε∂

x

ϕ

ε

= iξ(t)ϕ

ε

+ O √ ε e

Ct

in L

2

( R ), we write,

L e

ε

= i ∂

t

χ

1

+ ξ(t)∂

x

χ

1

ϕ

ε

+ O √ ε e

Ct

in L

2

(R).

The choice of the time-dependent eigenvectors ensures that for all time, the O (1) contribution of L e

ε

is orthogonal to the first mode (the eigenspace associated with λ

1

). Then, to get rid of these terms, we introduce a correction term to w

ε

. We set

θ

ε

(t, x) = w

ε

(t, x) + εg

ε

(t, x), g

ε

(t, x) = X

26j6P

X

16ℓ6dj

g

εj,ℓ

(t, x)χ

j

(x), where for j > 2 and for 1 6 ℓ 6 d

j

, the function g

j,ℓε

(t, x) solves the scalar Schr¨odinger equation

(1.16) iε∂

t

g

εj,ℓ

+ ε

2

2 ∂

x2

g

j,ℓε

− λ

j

(x)g

εj,ℓ

= ϕ

ε

r

j,ℓ

; g

εj,ℓ|t=0

= 0, where

(1.17) r

j,ℓ

(t, x) = − i ∂

t

χ

1

(t, x) + ξ(t)∂

x

χ

1

(t, x) , χ

j

(x)

CN

.

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The function θ

ε

(t) then solves

(1.18)

 

iε∂

t

θ

ε

(t, x) + ε

2

2 ∂

x2

θ

ε

(t, x) = V (x)θ

ε

(t, x) + εN L

ε

(t, x) + εL

ε

(t, x), θ

ε|t=0

= r

ε0

,

with

N L

ε

= Λ ε

1/2

| ϕ

ε

χ

1

+ θ

ε

− εg

ε

|

2CN

ε

χ

1

+ θ

ε

− εg

ε

) − | ϕ

ε

|

2

ϕ

ε

χ

1

, (1.19)

L

ε

= L e

ε

+

iε∂

t

+ ε

2

2 ∂

x2

− V (x)

g

ε

(t, x) (1.20)

= O ( √

εe

Ct

) + X

26j6P

X

16ℓ6dj

ε

2

2 ∂

x2

, χ

j

g

εj,ℓ

where the O ( √

εe

Ct

) holds in L

2

. The proof of the theorem then follows from a precise control of the functions χ

j

and g

εj,ℓ

, which is achieved in Sections 2 and 3, respectively. Then, the analysis of θ

ε

as ε goes to zero by an energy method is presented in Section 4.

2. The family of time-dependent eigenvectors

In this section we prove Proposition 1.8, recalling the construction of the eigen- vectors satisfying (1.14), and analyzing the behavior of their derivatives for large time. We follow the proof of [8]. More generally, we prove the following result which implies Proposition 1.8. We consider the Hamiltonian curves of

12

| ξ |

2

+ λ

j

(x), that we denote by (x

j

(t), ξ

j

(t)).

Proposition 2.1. There exists a smooth orthonormal basis of C

N

χ

j

(t, x)

16ℓ6dj

16j6P

such that for all t, χ

j

(t, x)

16ℓ6dj

spans the eigenspace associated to λ

j

, with χ

1

(0, x) = χ(x) and for m ∈ { 1, · · · , d

j

} ,

χ

mj

(t, x), ∂

t

χ

j

(t, x) + ξ

j

(t)∂

x

χ

j

(t, x)

CN

= 0.

Moreover, for ℓ ∈ { 1, . . . , d

j

} , k, p ∈ N, there exists a constant C = C(p, k) such that

tp

kx

χ

j

(t, x)

CN

6 C e

Ct

h x i

(k+p)(1+n0)

,

where n

0

appears in (1.2).

Proof of Proposition 2.1. We consider a smooth basis of eigenvectors (χ

j

(0))

16ℓ6dj

16j6P

such that χ

11

(0) = χ. Then, we denote by Π

j

(x) the smooth eigenprojector associ- ated with the eigenvalue λ

j

(x) and define

K

j

(x) = − i [Π

j

(x), ∂

x

Π

j

(x)] .

We set z = x − x

j

(t) and we consider the Schr¨odinger type equation

(2.1) i∂

t

Y

j

(t, z) = ξ

j

(t)K

j

(z + x

j

(t))Y

j

(t, z) ; Y

j

(0, z) = χ

j

(x

j

(0) + z).

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Let us prove that the vector Y

j

(t, z) is in the eigenspace of λ

j

(x

j

(t) + z). Indeed, the evolution of Z

j

(t, z) = (Id − Π

j

(x

j

(t) + z)) Y

j

(t, z) obeys to Z

j

(0, z) = 0 and

t

Z

j

(t, z) = − ξ

j

(t)∂

x

Π

j

(x

j

(t) + z)Y

j

− ξ

j

(t)(Id − Π

j

(x

j

(t) + z))[Π

j

(x

j

(t) + z), ∂

x

Π

j

(x

j

(t) + z)]Y

j

= − ξ

j

(t)∂

x

Π

j

(x

j

(t) + z)(Id − Π

j

(x

j

(t) + z))Y

j

= − ξ

j

(t)∂

x

Π

j

(x

j

(t) + z)Z

j

where we have used ∂

x

Π

j

= ∂

x

2j

) = Π

j

x

Π

j

+ (∂

x

Π

j

) Π

j

, whence Π

j

(∂

x

Π

j

j

= Π

j

j

x

Π

j

+ (∂

x

Π

j

) Π

j

) Π

j

= 2Π

j

(∂

x

Π

j

j

= 0.

Therefore, Z

j

(t) satisfies an equation of the form ∂

t

Z

j

= A(t, z)Z

j

, which combined with Z

j

(0) = 0, implies Z

j

(t) = 0 for all t ∈ R: the vectors Y

j

(t, z) are eigenvectors of V (x

j

(t) + z) for the eigenvalue λ

j

(x

j

(t) + z).

Besides, since ξ

j

(t)K

j

(z + x

j

(t)) is self-adjoint, Y

j

(t, z) is normalized for all t, and the family (Y

j

)

16ℓ6dj

is orthonormal. We define χ

j

(t, x) by

(2.2) χ

j

(t, x) = Y

j

(t, x − x

j

(t)) and we obtain an orthonormal basis of eigenvectors of V (x).

It remains to check that (1.14) holds. We have

t

χ

j

+ ξ

j

(t)∂

x

χ

j

= ∂

t

Y

j

(t, x − x

j

(t))

= iξ

j

(t)K

j

(x)χ

j

= − ξ

j

(t)[Π

j

(x), ∂

x

Π

j

(x)]χ

j

, whence

t

χ

j

+ ξ

j

(t)∂

x

χ

j

, χ

kj

CN

= − ξ

j

(t) [Π

j

, ∂

x

Π

j

j

, χ

kj

CN

= − ξ

j

(t) Π

j

j

, ∂

x

Π

j

j

χ

j

, χ

kj

CN

since χ

ℓ/kj

= Π

j

χ

ℓ/kj

. We then observe that Π

2j

= Π

j

implies Π

j

j

, ∂

x

Π

j

j

= Π

2j

x

Π

j

Π

j

− Π

j

x

Π

j

Π

2j

= 0.

This concludes the first part of Proposition 1.8. It remains to study the behavior at infinity of the vectors χ

j

(t, x) and of their derivatives.

By the definition of χ

j

(t, x) in (2.2), it is enough to prove

| ∂

tp

xk

Y

j

(t, z) |

CN

. e

Ct

h x

j

(t) + z i

(p+k)(1+n0)

.

For this, we crucially use the estimates of Lemma C.2 and we argue by induction.

Let us first consider the case p = 1 and k = 0. By Lemma C.2, we have | K

j

(x) | . h x i

1+n0

, whence (2.1) gives

| ∂

t

Y

j

(t, x − x

j

(t)) | . | ξ

j

(t) || K

j

(x) | . e

Ct

h x i

1+n0

. Let us now suppose k > 1 and p = 0. We observe that ∂

zk

Y

j

(t, z) solves (2.3)

( i∂

t

zk

Y

j

(t, z) = − iξ

j

(t)K

j

(z + x

j

(t))∂

zk

Y

j

(t, z) + f (t, z),

zk

Y

j

(0, z) = ∂

xk

χ

j

(0, z + x

j

(0)),

(10)

where

f (t, z) = X

06γ6k−1

c

γ

ξ

j

(t)∂

zγ

K

j

(z + x

j

(t))∂

zk−γ

Y

j

(t, z), for some complex numbers c

γ

independent of t and z. We obtain

zk

Y

j

(t, z) = U

j

(t, 0)∂

xk

χ

j

(z + x

j

(0)) + Z

t

0

U

j

(t, s)f (s, z)ds,

where U

j

(t, s) denotes the unitary propagator associated to (2.1) (when the initial time is equal to s). We have by Lemma C.2

| ∂

zk

Y

j

(0, z) |

CN

= ∂

xk

χ

j

(0, z + x

j

(0))

CN

. h z + x

j

(0) i

k(1+n0)

, therefore the induction assumption

∀ γ ∈ { 0, . . . , k − 1 } , | ∂

xγ

Y

j

(t, z) |

CN

. e

Ct

h x

j

(t) + z i

γ(1+n0)

implies, along with Lemma C.2,

| ∂

xk

Y

j

(t, z) |

CN

. e

Ct

h x

j

(t) + z i

k(1+n0)

.

We have obtained the estimate for p = 0, k ∈ N, and for p = 1, k = 0. Note that Equation (2.3) yields

∀ k ∈ N, | ∂

t

xk

Y

j

(t, z) |

CN

. e

Ct

h x

j

(t) + z i

(1+k)(1+n0)

,

and allows to prove the general estimate for time derivatives by an induction ar- gument which crucially uses the fact that we have an exponential control of the derivatives in time of ξ

j

(t). This property follows by induction from (1.7), (1.8),

and the fact that λ

j

is at most quadratic.

Before concluding this section, note that in view of the definition of the function r

j,ℓ

in (1.17), Proposition 1.8 gives the following corollary.

Corollary 2.2. For all p ∈ N and k ∈ N, there exists a constant C = C(p, k) such that, for x ∈ R, j ∈ { 1, · · · , P } and ℓ ∈ { 1, · · · , d

j

} ,

tp

kx

r

j,ℓ

(t, x) . e

Ct

h x i

(1+p+k)(1+n0)

. 3. Analysis of the correction terms

In this section, we will make use of the following norms defined for p ∈ N, k f k

Σpε

= sup

α+β6p

| x |

α

ε

β

f

(β)

(x)

L2

. We associate with this norm the functional space Σ

pε

defined by

Σ

pε

= { f ∈ L

2

(R

d

), k f k

Σpε

< ∞} .

In view of (1.8) and (1.13), for all p ∈ N, there exists c(p) such that (3.1) k ϕ

ε

(t) k

Σpε

. e

c(p)t

, ∀ t > 0.

We can obviously take c(0) = 0 by conservation of the L

2

-norm, but in general, the norm of ϕ

ε

in Σ

1ε

potentially grows exponentially in time (see [3]). We denote by U

kε

(t) the semi-group associated with the operator −

ε22

x2

+ λ

k

(x) and we observe that for p ∈ N, there exists a constant C(p) such that

(3.2) k U

kε

(t) k

L(Σpε)

6 C(p)e

C(p)|t|

.

The following averaging lemma shows an asymptotic orthogonality property.

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Lemma 3.1. For T > 0 and k 6 = j, there exists a constant C such that

∀ t ∈ [0, T ], ∀ p ∈ N , 1 iε

Z

t 0

U

kε

( − s)U

jε

(s)ds

L

Σ(p+3)nε 0 +p+2pε

6 Ce

Ct

.

Proof. We first observe that (3.3) iε∂

t

U

kε

( − t)U

jε

(t)

= U

kε

( − t) (λ

j

(x) − λ

k

(x)) U

jε

(t).

Indeed, if f ∈ L

2

(R) and f

ε

(t) = U

kε

( − t)U

jε

(t)f. We have iε∂

t

f

ε

(t, x) = −

− ε

2

2 ∂

x2

+ λ

k

(x)

f

ε

(t) + U

kε

( − t)

− ε

2

2 ∂

x2

+ λ

j

(x)

U

jε

(t)f

= U

kε

( − t) (λ

j

(x) − λ

k

(x)) U

jε

(t)f

because U

kε

( − t) commutes with −

ε22

x2

+ λ

k

(x). We use Equation (3.3) to perform an integration by parts:

U

kε

( − t)U

jε

(t) = U

kε

( − t) (λ

j

− λ

k

)

−1

U

kε

(t)U

kε

( − t) (λ

j

− λ

k

) U

jε

(t)

= iε U

kε

( − t) (λ

j

− λ

k

)

−1

U

kε

(t) ∂

t

U

kε

( − t)U

jε

(t) . Therefore,

1 iε

Z

t 0

U

kε

( − s)U

jε

(s)ds = h

U

kε

( − s) (λ

j

− λ

k

)

−1

U

jε

(s) i

t 0

− Z

t

0

s

U

kε

( − s) (λ

j

− λ

k

)

−1

U

kε

(s)

U

kε

( − s)U

jε

(s) ds.

Set

(3.4) γ

j,k

= (λ

k

− λ

j

)

−1

.

The behavior as x goes to infinity of these functions is studied in Appendix C (see Lemma C.1). It is proven there that for all β ∈ N,

xβ

γ

j,k

(x) . h x i

n0+|β|(1+n0)

.

Since the propagators U

kε

(t) and U

jε

(t) map continuously Σ

pε

into itself uniformly with respect to ε, we have

U

kε

( − s)γ

j,k

U

jε

(s)

t 0

L

Σ(p+1)nε 0 +ppε

. C(p),

where in all this paragraph, C(p) denotes a generic constant depending only on the parameter p ∈ N. Besides, we observe that

s

(U

kε

( − s)γ

j,k

U

kε

(s)) = 1

iε U

kε

( − s)

− ε

2

2 ∂

x2

+ λ

k

, γ

j,k

U

kε

(s).

In view of 1 iε

− ε

2

2 ∂

x2

+ λ

k

, γ

j,k

= 1 iε

− ε

2

2 ∂

x2

, γ

j,k

= iγ

j,k

(x)ε∂

x

+ iεγ

j,k′′

(x), and of

U

kε

( − s)γ

j,k

(x)ε∂

x

U

jε

(s)

L

Σ(p+3)nε 0 +p+2pε

+ U

kε

( − s)γ

′′j,k

(x)U

jε

(s)

L

Σ(p+2)nε 0 +p+2pε

. e

Cs

,

(12)

which comes from (3.2) and Lemma C.1, we get

k ∂

s

(U

kε

( − s)γ

j,k

U

kε

(s)) k

L(Σp+2εpε)

. e

Ct

,

which concludes the proof.

We now prove the following proposition.

Proposition 3.2. For p ∈ N , there exists C(p) such that for all j > 2, and all ℓ ∈ { 1, . . . , d

j

} ,

k g

j,ℓε

(t) k

Σpε

. e

C(p)t

, ∀ t > 0, where g

j,ℓε

is defined in (1.16).

Proof. We use Duhamel’s formula and write g

j,ℓε

(t) = 1

iε Z

t

0

U

jε

(t − s) (ϕ

ε

(s)r

j,ℓ

(s)) ds.

Besides, if ˜ ϕ

εj,ℓ

(t, x) = ϕ

ε

(t, x)r

j,ℓ

(t, x), then we have,

iε∂

t

+ ε

2

2 ∂

x2

− λ

1

(x)

˜

ϕ

εj,ℓ

= iε∂

t

r

j,ℓ

ϕ

ε

+ r

j,ℓ

ε

3/2

| ϕ

ε

|

2

ϕ

ε

+ ε

2

2

x2

, r

j,ℓ

(t, x) ϕ

ε

| {z }

=:εerε(t,x)

.

Therefore, we can write

˜

ϕ

εj,ℓ

(t) = U

1ε

(t) ˜ ϕ

εj,ℓ

(0) − i Z

t

0

U

1ε

(t − s) e r

ε

(s)ds, whence

g

j,ℓε

(t) = 1 iε

Z

t 0

U

jε

(t − s)U

1ε

(s)ds ϕ ˜

εj,ℓ

(0) − 1 ε

Z

t 0

Z

s 0

U

jε

(t − s)U

1ε

(s − τ) r e

ε

(τ)dτ ds

= 1 iε

Z

t 0

U

jε

(t − s)U

1ε

(s)ds ϕ ˜

εj,ℓ

(0) − Z

t

0

1 ε

Z

t τ

U

jε

(t − s)U

1ε

(s − τ)ds

e r

ε

(τ)dτ.

Lemma 3.1 yields

k g

j,ℓε

(t) k

Σpε

. e

Ct

+ Z

t

0

e

ke r

ε

(τ) k

Σqε

dτ,

with q = p + 2 + (p + 3)(1 + n

0

). Let us now study e r

ε

. We write e r

ε

= e r

ε1

+ e r

2ε

with e

r

ε1

(t, x) = i∂

t

r

j,ℓ

ϕ

ε

+ ε 2

x2

, r

j,ℓ

(t, x) ϕ

ε

. In view of Corollary 2.2 and of (3.1), we have for all q ∈ N ,

ke r

ε1

(t) k

Σqε(R)

. e

C(q)t

. A very rough estimate yields

ke r

2ε

(t) k

Σqε

= k √

εr

j,ℓ

| ϕ

ε

|

2

ϕ

ε

k

Σqε

. √

ε k r

j,ℓ

h x i

q

ϕ

ε

k

Σqε

kh ε∂

x

i

q

ϕ

ε

k

2L

. √

εe

Ct

h x i

q+(1+q)(n0+1)

ϕ

ε

Σqε

kh ε∂

x

i

q

ϕ

ε

k

2L

,

where we have used Corollary 2.2. Now with (1.13) and (3.1), we conclude ke r

2ε

(t) k

Σqε

. e

Ct

.

This completes the proof of Proposition 3.2.

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4. Consistency

We now prove Theorem 1.9. We go back to Equation (1.18), that we recall:

 

iε∂

t

θ

ε

(t, x) + ε

2

2 ∂

x2

θ

ε

(t, x) = V (x)θ

ε

(t, x) + εN L

ε

(t, x) + εL

ε

(t, x), θ

|t=0ε

= r

ε0

,

where N L

ε

and L

ε

are defined in (1.19) and (1.20), respectively. The standard L

2

-estimate yields:

k θ

ε

(t) k

L2

6 k r

ε0

k

L2

+ Z

t

0

( k N L

ε

(s) k

L2

+ k L

ε

(s) k

L2

) ds.

In view of (1.20), Proposition 2.1 and Proposition 3.2, we have k L

ε

(t) k

L2

. √

εe

Ct

. Besides, we observe

k N L

ε

(t) k

L2

. √

ε | ϕ

ε

(t) |

2

+ | θ

ε

(t) |

2CN

+ ε

2

| g

ε

(t) |

2CN

ε

(t) − εg

ε

(t))

L2

. √

ε k ϕ

ε

(t) k

2L

+ k θ

ε

(t) k

2L

+ ε

2

k g

ε

(t) k

2L

( k θ

ε

(t) k

L2

+ ε k g

ε

(t) k

L2

) . In view of (1.13), we have k ϕ

ε

(t) k

L

. ε

−1/4

e

Ct

. On the other hand, Proposi- tion 3.2 implies, in view of the Gagliardo-Nirenberg inequality

(4.1) k f k

L

. ε

−1/2

k f k

1/2L2

k ε∂

x

f k

1/2L2

, the estimate

ε

2

k g

ε

(t) k

2L

. εe

Ct

.

Therefore, it is natural to perform a bootstrap argument assuming, say (4.2) k θ

ε

(t) k

L

6 ε

−1/4

e

Ct

.

Note that we fixed the value of the constant in factor of the right hand side equal to one. We did so because θ

ε

, as an error term, is expected to be smaller than ϕ

ε

(the approximate solution) in the limit ε → 0. As long as (4.2) holds, the L

2

-estimate implies, in view of (1.5)

k θ

ε

(t) k

L2

. ε

κ

+ Z

t

0

√ εe

Cs

+ e

Cs

k θ

ε

(s) k

L2

ds.

By Gronwall Lemma, we obtain

(4.3) k θ

ε

(t) k

L2

6 C ε

κ

+ √ ε

e

eCt

.

It remains to check how long the bootstrap assumption (4.2) holds. For this, we use Gagliardo-Nirenberg inequality (4.1), and we look for a control of the norm of θ

ε

(t) in Σ

1ε

. Differentiating the system (1.18) with respect to x, we find

iε∂

t

(ε∂

x

θ

ε

) + ε

2

2 ∂

x2

(ε∂

x

θ

ε

) = V (x)ε∂

x

θ

ε

+ εV

(x)θ

ε

+ ε

2

x

N L

ε

+ ε

2

x

L

ε

, We observe that since V is at most quadratic, | V

(x)θ

ε

|

CN

. h x i | θ

ε

|

CN

. Therefore, in order to obtain a closed system of estimates, we consider the equation satisfied by xθ

ε

: multiply (1.18) by x,

iε∂

t

(xθ

ε

) + ε

2

2 ∂

x2

(xθ

ε

) = V (x)(xθ

ε

) + ε

2

x

θ

ε

+ εxN L

ε

+ εxL

ε

.

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By Proposition 3.2, we have

k xL

ε

(t) k

L2

+ k ε∂

x

L

ε

(t) k

L2

. √ εe

Ct

. Besides,

| xN L

ε

(t, x) |

CN

. | φ

ε

(t, x) |

2

+ | θ

ε

(t, x) |

2CN

+ ε

2

| g

ε

(t, x) |

2CN

×

× ( | xθ

ε

(t, x) |

CN

+ ε | xg

ε

(t, x) |

CN

) ,

| ε∂

x

N L

ε

(t, x) |

CN

. | φ

ε

(t, x) |

2

+ | θ

ε

(t, x) |

2CN

+ ε

2

| g

ε

(t, x) |

2CN

| ε∂

x

θ

ε

(t, x) |

CN

+ ε | φ

ε

(t, x) |

2

+ | θ

ε

(t, x) |

2CN

+ ε

2

| g

ε

(t, x) |

2CN

| ε∂

x

g

ε

(t, x) |

CN

+ | ε∂

x

φ

ε

(t, x) | × | φ

ε

(t, x) | × | θ

ε

(t, x) |

CN

+ ε | φ

ε

(t, x) |

2

× | ∂

x

χ

1

(t, x) |

CN

× | θ

ε

(t, x) |

CN

.

Arguing as before and using again (1.13), we obtain that under (1.2) we have k ε∂

x

θ

ε

(t) k

L2

+ k xθ

ε

(t) k

L2

. ε

κ

+ √

ε e

eCt

. Gagliardo–Nirenberg inequality then implies

k θ

ε

(t) k

L

. ε

−1/2

ε

κ

+ √ ε

e

eCt

. We infer that (4.2) holds (at least) as long as

ε

κ−1/2

+ 1

e

eCt

≪ ε

−1/4

e

Ct

, which is ensured provided that t 6 Cloglog

1ε

, for some suitable constant C, since κ > 1/4. This concludes the bootstrap argument: we infer

sup

|t|6Cloglog

(

1ε

)

( k θ

ε

(t) k

L2

+ k xθ

ε

(t) k

L2

+ k ε∂

x

θ

ε

(t) k

L2

) −→

ε→0

0.

Theorem 1.9 then follows from the above asymptotics, together with the relation θ

ε

= w

ε

+ εg

ε

, and Proposition 3.2.

Remark 4.1. In the case where V

ε

= D + εW , as in Remark 1.10, the proof can be adapted, in order to reproduce the argument given in [4]. The main point to notice is that (local in time) Strichartz estimates are available for the propagator associated to −

ε22

x2

+ D(x), thanks to [6]. Then in the presence of the power ε in front of W , the potential εW can be considered as a source term in the error estimates: the factor ε is crucial to avoid a singular power of ε due to the presence of ε in front of the time derivative in (1.18). The proof in [4, Section 6] for the cubic, one-dimensional Schr¨odinger equation can be reproduced: another bootstrap argument can be invoked, which does not involve Gagliardo–Nirenberg inequalities, since a useful a priori estimate for the envelope u is available.

5. Superposition

As explained in the introduction, the only difficulty in the proof of Theorem 1.11 is to treat a nonlinear interaction term. Indeed, we set

w

ε

= ψ

ε

− ϕ

ε1

χ

1

− ϕ

ε2

χ

2

+ εg

ε

where g

ε

is the sum of two correction terms, similar to the one introduced in § 1.3.

More precisely, set p(1) = 1, and p(2) = 1 if ˜ λ

1

= ˜ λ

2

, p(2) = 2 otherwise. Define

(15)

g

ε

= g

1ε

+ g

ε2

, with

g

1ε

= X

16j6P, j6=p(1)

X

16ℓ6dj

g

j,1,ℓε

(t, x)χ

j

(t, x), g

2ε

= X

16j6P, j6=p(2)

X

16ℓ6dj

g

j,2,ℓε

(t, x)χ

j

(t, x),

where for k = { 1, 2 } , j 6 = p(k) and 1 6 ℓ 6 d

j

, the function g

εj,k,ℓ

(t, x) solves the scalar Schr¨odinger equation

(5.1) iε∂

t

g

j,k,ℓε

+ ε

2

2 ∂

x2

g

εj,k,ℓ

− λ

j

(x)g

εj,k,ℓ

= ϕ

ε

r

j,k,ℓ

; g

j,k,ℓ|t=0ε

= 0, where

(5.2) r

j,k,ℓ

(t, x) = − i ∂

t

χ

k

(t, x) + ξ

p(k)

(t)∂

x

χ

k

(t, x) , χ

j

(t, x)

CN

. The function w

ε

(t) then solves

iε∂

t

w

ε

+ ε

2

2 ∂

x2

w

ε

= V (x)w

ε

+ εN L

ε

+ εL

ε

; w

ε|t=0

= 0, with

L

ε

= O ( √

εe

Ct

) + X

k=1,2

X

16j6P j6=p(k)

X

16ℓ6dj

ε

2

2 ∂

x2

, χ

j

g

j,k,ℓε

= O ( √ εe

Ct

).

Here, the O ( √

εe

Ct

) holds in Σ

1ε

, from Proposition 3.2. Besides, N L

ε

= √

ε w

ε

+ ϕ

ε1

χ

1

+ ϕ

ε2

χ

2

+ εg

ε

2

w

ε

+ ϕ

ε1

χ

1

+ ϕ

ε2

χ

2

+ εg

ε

− | ϕ

ε1

|

2

ϕ

ε1

χ

1

− | ϕ

ε2

|

2

ϕ

ε2

χ

2

Adding and subtracting the term √

ε | ϕ

ε1

χ

1

+ ϕ

ε2

χ

2

|

2

ε1

χ

1

+ ϕ

ε2

χ

2

), we have

| N L

ε

| 6 N

Sε

+ N

Iε

, where we have the pointwise estimates

N

Iε

. √

ε | ϕ

ε1

|

2

| ϕ

ε2

| + | ϕ

ε2

|

2

| ϕ

ε1

| , N

Sε

. √

ε | ϕ

ε1

|

2

+ | ϕ

ε2

|

2

+ | w

ε

|

2

+ ε

2

| g

ε

|

2

( | w

ε

| + ε | g

ε

| ) .

The semilinear term N

Sε

can be treated exactly in the same manner as in Section 4.

It remains to analyze Z

t

0

k N L

εI

(s) k

Σ1ε

ds. We observe

√ ε Z

t

0

| ϕ

ε1

(s) |

2

ϕ

ε2

(s)

L2

ds = Z

t

0

u

1

s, y − x

1

(s) − x

2

(s)

√ ε

2

u

2

(s, y)

L2

ds,

and we note that the contribution of | ϕ

ε1

|

2

ϕ

ε2

and that of | ϕ

ε2

|

2

ϕ

ε1

play the same role. Also, we leave out the other terms which are needed in view of a Σ

1ε

estimate, since they create no trouble. Arguing as in [4, Lemma 6.1], we obtain:

Lemma 5.1. Let T ∈ R, 0 < γ < 1/2 and

I

ε

(T ) = { t ∈ [0, T ], | x

1

(t) − x

2

(t) | 6 ε

γ

} .

(16)

Then, for all k ∈ N, there exists a constant C

k

such that Z

T

0

k N L

εI

(t) k

Σ1ε

dt . (M

k+2

(T ))

3

T ε

k(1/2−γ)

+ | I

ε

(T ) | e

CkT

, with

M

k

(T ) = sup

k h x i

α

xβ

u

j

k

L([0,T],L2(R))

; j ∈ { 1, 2 } , α + β 6 k . In view of this lemma and of Equation (1.12), we obtain

Z

T

0

k N L

εI

(t) k

Σ1ε

dt . e

CT

T ε

k(1/2−γ)

+ | I

ε

(T ) | ,

and the next lemma yields the conclusion.

Lemma 5.2. Set

Γ = inf

x∈R

λ ˜

1

(x) − λ ˜

2

(x) − (E

1

− E

2

) ,

and suppose Γ > 0. Then for 0 < γ < 1/2, there exists C

0

, C

1

> 0 such that

| I

ε

(t) | . ε

γ

Γ

−2

e

C0t

, 0 6 t 6 C

1

log 1

ε

.

Proof. Consider J

ε

(t) an interval of maximal length included in I

ε

(t), and N

ε

(t) the number of such intervals. The result comes from the estimate

| I

ε

(t) | 6 N

ε

(t) × max | J

ε

(t) | , with

(5.3) | J

ε

(t) | . ε

γ

e

Ct

Γ

−1

and N

ε

(t) . te

Ct

Γ

−1

,

provided that ε

γ

e

Ct

≪ 1. Let us prove the first property: consider τ, σ ∈ J

ε

(t).

There exists t

∈ [τ, σ] such that

| (x

1

(τ) − x

2

(τ)) − (x

1

(σ) − x

2

(σ)) | = | τ − σ | | ξ

1

(t

) − ξ

2

(t

) | , whence

| τ − σ | 6 | ξ

1

(t

) − ξ

2

(t

) |

−1

× 2ε

γ

. On the other hand,

| ξ

1

(t

) − ξ

2

(t

) | > || ξ

1

(t

) | − | ξ

2

(t

) || >

| ξ

1

(t

) |

2

− | ξ

2

(t

) |

2

| ξ

1

(t

) | + | ξ

2

(t

) | . We use

| ξ

1

(t

) | + | ξ

2

(t

) | . e

Ct

,

| ξ

1

(t

) |

2

− | ξ

2

(t

) |

2

= 2

E

1

− E

2

− ˜ λ

1

(x

1

(t

)) + ˜ λ

2

(x

2

(t

)) , and infer

E

1

− E

2

− λ ˜

1

(x

1

(t

)) + ˜ λ

2

(x

2

(t

)) >

E

1

− E

2

− ˜ λ

1

(x

1

(t

)) + ˜ λ

2

(x

1

(t

))

− λ ˜

2

(x

1

(t

)) + ˜ λ

2

(x

2

(t

))

> Γ − Cε

γ

e

Ct

,

(17)

where we have used the fact that ˜ λ

2

is at most quadratic. Therefore, if ε

γ

e

Ct

is sufficiently small,

E

1

− E

2

− λ ˜

1

(x

1

(t

)) + ˜ λ

2

(x

2

(t

)) > Γ 2 . We infer

| τ − σ | . ε

γ

e

Ct

Γ

−1

, provided ε

γ

e

Ct

≪ 1.

Let us now consider N

ε

(t). We use that as t is large, N

ε

(t) is comparable to the number of distinct intervals of maximal size where | x

1

(t) − x

2

(t) | > ε

γ

. More precisely, N

ε

(t) is smaller than t divided by the minimal size of these intervals.

Therefore, we consider one interval ]τ, σ[ of this type and we look for lower bound of σ − τ. We have

| x

1

(τ) − x

2

(τ) | = | x

1

(σ) − x

2

(σ) | = ε

γ

, and ∀ t ∈ [τ, σ], | x

1

(t) − x

2

(t) | > ε

γ

. Besides, inside ]τ, σ[, x

1

(t) − x

2

(t) has a constant sign that we can suppose to be + (one argues similarly if it is − ). Under this assumption, we have

ξ

1

(τ) − ξ

2

(τ) > 0 and ξ

1

(σ) − ξ

2

(σ) < 0.

Using the exponential control of λ

j

(x

j

(t)) for j ∈ { 1, 2 } , we obtain (5.4) (ξ

1

(τ) − ξ

2

(τ)) − (ξ

1

(σ) − ξ

2

(σ)) . e

Ct

(σ − τ).

We write

ξ

1

(τ) − ξ

2

(τ) = | ξ

1

(τ) − ξ

2

(τ) | >

| ξ

1

(τ) |

2

− | ξ

2

(τ) |

2

| ξ

1

(τ) | + | ξ

2

(τ) | (5.5)

& e

−Ct

| ξ

1

(τ) |

2

− | ξ

2

(τ) |

2

and

(5.6) − ξ

1

(σ) + ξ

2

(σ) = | ξ

1

(τ) − ξ

2

(τ ) | & e

−Ct

| ξ

1

(σ) |

2

− | ξ

2

(σ) |

2

. As before, we prove

| ξ

1

(τ) |

2

− | ξ

2

(τ) |

2

+ | ξ

1

(σ) |

2

− | ξ

2

(σ) |

2

& Γ,

provided that ε

γ

e

Ct

≪ 1. Therefore, plugging the latter equation, (5.5) and (5.6) into (5.4), we obtain

σ − τ & e

−2Ct

Γ thus N

ε

(t) . te

Ct

Γ

−1

. e

Ct

Γ

−1

which completes the proof of Theorem 1.11.

Remark 5.3. The proof shows that if the approximation of Theorem 1.9 is proven

to be valid on some time interval [0, C log(1/ε)], then Theorem 1.11 will also be

valid on a time interval of the same form.

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